Share and discover more about AI with social posts from the community.huggingface/OpenAi
SAM 2: Segment Anything in Images and Videos
Segment Anything Model 2 (SAM 2) is a foundation model towards solving promptable visual segmentation in images and videos. We extend SAM to video by considering images as a video with a single frame. The model design is a simple transformer architecture with streaming memory for real-time video processing. We build a model-in-the-loop data engine, which improves model and data via user interaction, to collect our SA-V dataset, the largest video segmentation dataset to date. SAM 2 trained on our data provides strong performance across a wide range of tasks and visual domains.https://github.com/zsxkib/segment-anything-2/raw/video/assets/sa_v_dataset.jpg?raw=true
Segment Anything Model 2 (SAM 2) is a foundation model towards solving promptable visual segmentation in images and videos. We extend SAM to video by considering images as a video with a single frame. The model design is a simple transformer architecture with streaming memory for real-time video processing. We build a model-in-the-loop data engine, which improves model and data via user interaction, to collect our SA-V dataset, the largest video segmentation dataset to date. SAM 2 trained on our data provides strong performance across a wide range of tasks and visual domains.https://github.com/zsxkib/segment-anything-2/raw/video/assets/sa_v_dataset.jpg?raw=true
Batouresearch / high-resolution-controlnet-tile
Run time and cost
This model costs approximately $0.054 to run on Replicate, or 18 runs per $1, but this varies depending on your inputs. It is also open source and you can run it on your own computer with Docker.
This model runs on Nvidia A40 (Large) GPU hardware. Predictions typically complete within 75 seconds. The predict time for this model varies significantly based on the inputs.
Readme
High quality upscale from Fermat.app. Increase the creativity to encourage hallucination.https://replicate.delivery/pbxt/etT436Z2RrWAOajwhQm6YLBHiT5Y1Oix2aZnDLnIkfM7u4ESA/out-0.png
https://replicate.delivery/pbxt/1rbKAbFss7ZUGNxKnFGmOEHBaeEZ7cI7Sx61eiOo9AyGjQMTA/output.jpg
Run time and cost
This model costs approximately $0.054 to run on Replicate, or 18 runs per $1, but this varies depending on your inputs. It is also open source and you can run it on your own computer with Docker.
This model runs on Nvidia A40 (Large) GPU hardware. Predictions typically complete within 75 seconds. The predict time for this model varies significantly based on the inputs.
Readme
High quality upscale from Fermat.app. Increase the creativity to encourage hallucination.https://replicate.delivery/pbxt/etT436Z2RrWAOajwhQm6YLBHiT5Y1Oix2aZnDLnIkfM7u4ESA/out-0.png
https://replicate.delivery/pbxt/1rbKAbFss7ZUGNxKnFGmOEHBaeEZ7cI7Sx61eiOo9AyGjQMTA/output.jpg
Restore images
These models restore and improve images by fixing defects like blur, noise, and low resolution. Key capabilities:
Deblurring - Sharpen blurry images by reversing blur effects. Useful for old photos.
Denoising - Remove grain and artifacts by learning noise patterns.
Colorization - Add realistic color to black and white photos.
Face restoration - Improve the image quality of faces in old photos, or unrealistic AI generated faces.
Our Picks
Best restoration model: google-research/maxim
If you need to sharpen a blurry photo, or remove noise or compression artifacts, start with google-research/maxim. It has a total of 11 image restoration models baked-in that let you deblur, denoise, remove raindrops, and more. If you’re not getting the results you’re looking for, try megvii-research/nafnet which is similar but supports fewer restoration features.
Best colorization model: piddnad/ddcolor
The best model for adding color to black and white photos is piddnad/ddcolor, which was released in 2023. If you are looking for more saturated results try out arielreplicate/deoldify_image.
Best face restoration model: sczhou/codeformer
If you’re looking for a face restoration model, try starting with sczhou/codeformer. It produces more realistic faces than alternatives like tencentarc/gfpgan. If you aren’t getting the exact image improvements you want, we recommend exploring more modern upscaling models like batouresearch/magic-image-refiner.https://replicate.com/collections/image-restoration
These models restore and improve images by fixing defects like blur, noise, and low resolution. Key capabilities:
Deblurring - Sharpen blurry images by reversing blur effects. Useful for old photos.
Denoising - Remove grain and artifacts by learning noise patterns.
Colorization - Add realistic color to black and white photos.
Face restoration - Improve the image quality of faces in old photos, or unrealistic AI generated faces.
Our Picks
Best restoration model: google-research/maxim
If you need to sharpen a blurry photo, or remove noise or compression artifacts, start with google-research/maxim. It has a total of 11 image restoration models baked-in that let you deblur, denoise, remove raindrops, and more. If you’re not getting the results you’re looking for, try megvii-research/nafnet which is similar but supports fewer restoration features.
Best colorization model: piddnad/ddcolor
The best model for adding color to black and white photos is piddnad/ddcolor, which was released in 2023. If you are looking for more saturated results try out arielreplicate/deoldify_image.
Best face restoration model: sczhou/codeformer
If you’re looking for a face restoration model, try starting with sczhou/codeformer. It produces more realistic faces than alternatives like tencentarc/gfpgan. If you aren’t getting the exact image improvements you want, we recommend exploring more modern upscaling models like batouresearch/magic-image-refiner.https://replicate.com/collections/image-restoration
Using FLUX.1 Schnell for faster inference
You can use your FLUX.1 Dev LoRA with the smaller FLUX.1 Schnell model, to generate images faster and cheaper. Just change the model parameter from “dev” to “schnell” when you generate, and lower your number of steps to something small like 4.
Note that outputs will still be under the non-commercial license of FLUX.1 Dev.
Examples and use cases
Check out our examples gallery for inspiration. You can see how others have fine-tuned FLUX.1 to create different styles, characters, a never-ending parade of cute animals, and more.https://d31rfu1d3w8e4q.cloudfront.net/static/blog/fine-tune-flux/3-base.webp
You can use your FLUX.1 Dev LoRA with the smaller FLUX.1 Schnell model, to generate images faster and cheaper. Just change the model parameter from “dev” to “schnell” when you generate, and lower your number of steps to something small like 4.
Note that outputs will still be under the non-commercial license of FLUX.1 Dev.
Examples and use cases
Check out our examples gallery for inspiration. You can see how others have fine-tuned FLUX.1 to create different styles, characters, a never-ending parade of cute animals, and more.https://d31rfu1d3w8e4q.cloudfront.net/static/blog/fine-tune-flux/3-base.webp
How to fine-tune FLUX.1
Fine-tuning FLUX.1 on Replicate is a straightforward process that can be done either through the web interface or programmatically via the API. Let’s walk through both methods.
Prepare your training data
To start fine-tuning, you’ll need a collection of images that represent the concept you want to teach the model. These images should be diverse enough to cover different aspects of the concept. For example, if you’re fine-tuning on a specific character, include images in various settings, poses, and lighting. Here are some guidelines:
Use 12-20 images for best results
Use large images if possible
Use JPEG or PNG formats
(Optional) Create a corresponding .txt file for each image with the same name, containing the caption
Once you have your images (and optional captions), zip them up into a single file.
Start the training process
https://replicate.com/blog/fine-tune-flux
Fine-tuning FLUX.1 on Replicate is a straightforward process that can be done either through the web interface or programmatically via the API. Let’s walk through both methods.
Prepare your training data
To start fine-tuning, you’ll need a collection of images that represent the concept you want to teach the model. These images should be diverse enough to cover different aspects of the concept. For example, if you’re fine-tuning on a specific character, include images in various settings, poses, and lighting. Here are some guidelines:
Use 12-20 images for best results
Use large images if possible
Use JPEG or PNG formats
(Optional) Create a corresponding .txt file for each image with the same name, containing the caption
Once you have your images (and optional captions), zip them up into a single file.
Start the training process
https://replicate.com/blog/fine-tune-flux
Replicate
Fine-tune FLUX.1 with your own images
We've added fine-tuning (LoRA) support to FLUX.1 image generation models. You can train FLUX.1 on your own images with one line of code using Replicate's API.
What is fine-tuning FLUX.1 on Replicate?
These big image generation models, like FLUX.1 and Stable Diffusion, are trained on a bunch of images that have had noise added, and they learn the reverse function of “adding noise.” Amazingly, that turns out to be “creating images.”
How do they know which image to create? They build on transformer models, like CLIP and T5, which are themselves trained on tons of image-caption pairs. These are language-to-image encoders: they learn to map an image and its caption to the same shape in high-dimensional space. When you send them a text prompt, like “squirrel reading a newspaper in the park,” they can map that to patterns of pixels in a grid. To the encoder, the picture and the caption are the same thing.
The image generation process looks like this: take some input pixels, move them a little bit away from noise and toward the pattern created by your text input, and repeat until the correct number of steps is reached.
The fine-tuning process, in turn, takes each image/caption pair from your dataset and updates that internal mapping a little bit. You can teach the model anything this way, as long as it can be represented through image-caption pairs: characters, settings, mediums, styles, genres. In training the model will learn to associate your concept with a particular text string. Include this string in your prompt to activate that association.
For example, say you want to fine-tune the model on your comic book superhero. You’ll collect some images of your character as your dataset. A well-rounded batch: different settings, costumes, lighting, maybe even different art styles. That way the model understands that what it’s learning is this one person, not any of these other incidental details.
Pick a short, uncommon word or phrase as your trigger: something unique that won’t conflict with other concepts or fine-tunes. You might choose something like “bad 70s food” or “JELLOMOLD”. Train your model. Now, when you prompt “Establishing shot of bad 70s food at a party in San Francisco,” your model will call up your specific concept. Easy as that.
Could it be as easy as that? Yes, actually. We realized that we could use the Replicate platform to make fine-tuning as simple as uploading images. We can even do the captioning for you.
If you’re not familiar with Replicate, we make it easy to run AI as an API. You don’t have to go looking for a beefy GPU, you don’t have to deal with environments and containers, you don’t have to worry about scaling. You write normal code, with normal APIs, and pay only for what you use.
You can try this right now! It doesn’t take a lot of images. Check out our examples gallery to see the kinds of styles and characters people are creating.
Grab a few photos of your pet, or your favorite houseplant, and let’s get started.
These big image generation models, like FLUX.1 and Stable Diffusion, are trained on a bunch of images that have had noise added, and they learn the reverse function of “adding noise.” Amazingly, that turns out to be “creating images.”
How do they know which image to create? They build on transformer models, like CLIP and T5, which are themselves trained on tons of image-caption pairs. These are language-to-image encoders: they learn to map an image and its caption to the same shape in high-dimensional space. When you send them a text prompt, like “squirrel reading a newspaper in the park,” they can map that to patterns of pixels in a grid. To the encoder, the picture and the caption are the same thing.
The image generation process looks like this: take some input pixels, move them a little bit away from noise and toward the pattern created by your text input, and repeat until the correct number of steps is reached.
The fine-tuning process, in turn, takes each image/caption pair from your dataset and updates that internal mapping a little bit. You can teach the model anything this way, as long as it can be represented through image-caption pairs: characters, settings, mediums, styles, genres. In training the model will learn to associate your concept with a particular text string. Include this string in your prompt to activate that association.
For example, say you want to fine-tune the model on your comic book superhero. You’ll collect some images of your character as your dataset. A well-rounded batch: different settings, costumes, lighting, maybe even different art styles. That way the model understands that what it’s learning is this one person, not any of these other incidental details.
Pick a short, uncommon word or phrase as your trigger: something unique that won’t conflict with other concepts or fine-tunes. You might choose something like “bad 70s food” or “JELLOMOLD”. Train your model. Now, when you prompt “Establishing shot of bad 70s food at a party in San Francisco,” your model will call up your specific concept. Easy as that.
Could it be as easy as that? Yes, actually. We realized that we could use the Replicate platform to make fine-tuning as simple as uploading images. We can even do the captioning for you.
If you’re not familiar with Replicate, we make it easy to run AI as an API. You don’t have to go looking for a beefy GPU, you don’t have to deal with environments and containers, you don’t have to worry about scaling. You write normal code, with normal APIs, and pay only for what you use.
You can try this right now! It doesn’t take a lot of images. Check out our examples gallery to see the kinds of styles and characters people are creating.
Grab a few photos of your pet, or your favorite houseplant, and let’s get started.
Fine-tune FLUX.1 with your own images
We have fine-tuning for FLUX.1
Fine-tuning is now available on Replicate for the FLUX.1 [dev] image model. Here’s what that means, and how to do it.
FLUX.1 is a family of text-to-image models released by Black Forest Labs this summer. The FLUX.1 models set a new standard for open-weights image models: they can generate realistic hands, legible text, and even the strangely hard task of “funny memes.” You can now fine-tune your model on Replicate with the FLUX.1 Dev LoRA Trainer.
If you know what all that means and you’re ready to try it with your dataset, you can skip to the code.
Otherwise, here’s what it means and why you should care.
https://replicate.com/blog/fine-tune-flux
We have fine-tuning for FLUX.1
Fine-tuning is now available on Replicate for the FLUX.1 [dev] image model. Here’s what that means, and how to do it.
FLUX.1 is a family of text-to-image models released by Black Forest Labs this summer. The FLUX.1 models set a new standard for open-weights image models: they can generate realistic hands, legible text, and even the strangely hard task of “funny memes.” You can now fine-tune your model on Replicate with the FLUX.1 Dev LoRA Trainer.
If you know what all that means and you’re ready to try it with your dataset, you can skip to the code.
Otherwise, here’s what it means and why you should care.
https://replicate.com/blog/fine-tune-flux
Replicate
Fine-tune FLUX.1 with your own images
We've added fine-tuning (LoRA) support to FLUX.1 image generation models. You can train FLUX.1 on your own images with one line of code using Replicate's API.
Welcome to our interactive tutorial on Tony's Chocolonely! I'm excited to guide you through the world of this Dutch chocolate company, known for its commitment to selling 100% slavery-free chocolate.
Overview of Tony's Chocolonely
Tony's Chocolonely is a Dutch chocolate company founded in 2005 by Teun van de Keuken, a Dutch TV journalist who was shocked to discover that most of the chocolate on supermarket shelves was harvested by slaves. He decided to take action and make 5,000 Fairtrade chocolate bars himself, marking the birth of Tony's Chocolonely.
Key Aspects of Tony's Chocolonely
Fairtrade Certification: Tony's Chocolonely is committed to Fairtrade certification, which guarantees a minimum price for cocoa farmers and promotes fair labor practices.
Slavery-Free Chocolate: Tony's Chocolonely aims to eradicate slavery from the global chocolate industry, making it a pioneer in the movement towards slavery-free chocolate.
Direct Trade: The company buys cocoa mass directly from cocoa farmers in West Africa, working with 5 cocoa cooperatives in Ghana and Ivory Coast.
Environmental Sustainability: Tony's Chocolonely prioritizes environmental sustainability, using uncoated, recycled FSC-certified paper for its wrappers and offsetting emissions through Justdiggit.
What would you like to learn more about?
Please choose one of the following topics:
History of Tony's Chocolonely: Learn about the company's founding, its mission, and its journey towards becoming a leader in the chocolate industry.
Fairtrade Certification and Slavery-Free Chocolate: Dive deeper into the Fairtrade certification process, the benefits of slavery-free chocolate, and Tony's Chocolonely's commitment to these values.
Direct Trade and Environmental Sustainability: Explore the company's direct trade practices, its environmental sustainability initiatives, and how these efforts contribute to a more responsible chocolate industry.
Customized Wrappers and Gift Cards: Discover how to create customized wrappers, order gift cards, and other fun features of Tony's Chocolonely's online shop.
Other: Is there something specific you'd like to know about Tony's Chocolonely?
Overview of Tony's Chocolonely
Tony's Chocolonely is a Dutch chocolate company founded in 2005 by Teun van de Keuken, a Dutch TV journalist who was shocked to discover that most of the chocolate on supermarket shelves was harvested by slaves. He decided to take action and make 5,000 Fairtrade chocolate bars himself, marking the birth of Tony's Chocolonely.
Key Aspects of Tony's Chocolonely
Fairtrade Certification: Tony's Chocolonely is committed to Fairtrade certification, which guarantees a minimum price for cocoa farmers and promotes fair labor practices.
Slavery-Free Chocolate: Tony's Chocolonely aims to eradicate slavery from the global chocolate industry, making it a pioneer in the movement towards slavery-free chocolate.
Direct Trade: The company buys cocoa mass directly from cocoa farmers in West Africa, working with 5 cocoa cooperatives in Ghana and Ivory Coast.
Environmental Sustainability: Tony's Chocolonely prioritizes environmental sustainability, using uncoated, recycled FSC-certified paper for its wrappers and offsetting emissions through Justdiggit.
What would you like to learn more about?
Please choose one of the following topics:
History of Tony's Chocolonely: Learn about the company's founding, its mission, and its journey towards becoming a leader in the chocolate industry.
Fairtrade Certification and Slavery-Free Chocolate: Dive deeper into the Fairtrade certification process, the benefits of slavery-free chocolate, and Tony's Chocolonely's commitment to these values.
Direct Trade and Environmental Sustainability: Explore the company's direct trade practices, its environmental sustainability initiatives, and how these efforts contribute to a more responsible chocolate industry.
Customized Wrappers and Gift Cards: Discover how to create customized wrappers, order gift cards, and other fun features of Tony's Chocolonely's online shop.
Other: Is there something specific you'd like to know about Tony's Chocolonely?
Welcome to our interactive tutorial on LoRA (Lightweight Adapter) and Diffusion Pipelines!
In this tutorial, we'll be exploring the AlbertFLUX model, a LoRA derived from the FLUX.1-dev model by black-forest-labs. We'll delve into the details of the model's architecture, training settings, and inference pipeline.
What would you like to learn about?
LoRA and its applications: Understand the concept of LoRA and its use cases in machine learning.
AlbertFLUX model architecture: Learn about the model's architecture, including the text encoder and image decoder.
Training settings and hyperparameters: Explore the training settings and hyperparameters used to train the AlbertFLUX model.
Inference pipeline and code: Understand the inference pipeline and code used to generate images with the AlbertFLUX model.
Diffusion pipelines and their benefits: Learn about diffusion pipelines and their advantages in image generation tasks.
Please select one of the above options to get started!
https://huggingface.co/markury/AlbertFLUX
In this tutorial, we'll be exploring the AlbertFLUX model, a LoRA derived from the FLUX.1-dev model by black-forest-labs. We'll delve into the details of the model's architecture, training settings, and inference pipeline.
What would you like to learn about?
LoRA and its applications: Understand the concept of LoRA and its use cases in machine learning.
AlbertFLUX model architecture: Learn about the model's architecture, including the text encoder and image decoder.
Training settings and hyperparameters: Explore the training settings and hyperparameters used to train the AlbertFLUX model.
Inference pipeline and code: Understand the inference pipeline and code used to generate images with the AlbertFLUX model.
Diffusion pipelines and their benefits: Learn about diffusion pipelines and their advantages in image generation tasks.
Please select one of the above options to get started!
https://huggingface.co/markury/AlbertFLUX
huggingface.co
markury/AlbertFLUX · Hugging Face
We’re on a journey to advance and democratize artificial intelligence through open source and open science.
Welcome to our interactive tutorial on Flux DreamBooth LoRA!
We will be exploring the concept of LoRA (Low-Rank Adaptation) in the context of DreamBooth, a text-to-image synthesis model. Specifically, we will be working with the linoyts/3d_icon_flux_1500 DreamBooth LoRA weights for the black-forest-labs/FLUX.1-dev model.
What would you like to learn about?
Model Description: Understand the context and training process of the Flux DreamBooth LoRA model.
LoRA and its Applications: Learn about the concept of LoRA and its use cases in text-to-image synthesis.
Trigger Words and Image Generation: Discover how to use trigger words to generate images with the Flux DreamBooth LoRA model.
Loading and Using the Model: Understand how to load and use the linoyts/3d_icon_flux_1500 DreamBooth LoRA weights with the diffusers library.
Limitations and Bias: Explore the potential limitations and biases of the model, as well as strategies for addressing them.
Training Details: Learn about the data used to train the model and its implications for the generated images.
Please select one of the above options by typing the corresponding number.
https://huggingface.co/linoyts/3d_icon_flux_1500#flux-dreambooth-lora---linoyts3d_icon_flux_1500
We will be exploring the concept of LoRA (Low-Rank Adaptation) in the context of DreamBooth, a text-to-image synthesis model. Specifically, we will be working with the linoyts/3d_icon_flux_1500 DreamBooth LoRA weights for the black-forest-labs/FLUX.1-dev model.
What would you like to learn about?
Model Description: Understand the context and training process of the Flux DreamBooth LoRA model.
LoRA and its Applications: Learn about the concept of LoRA and its use cases in text-to-image synthesis.
Trigger Words and Image Generation: Discover how to use trigger words to generate images with the Flux DreamBooth LoRA model.
Loading and Using the Model: Understand how to load and use the linoyts/3d_icon_flux_1500 DreamBooth LoRA weights with the diffusers library.
Limitations and Bias: Explore the potential limitations and biases of the model, as well as strategies for addressing them.
Training Details: Learn about the data used to train the model and its implications for the generated images.
Please select one of the above options by typing the corresponding number.
https://huggingface.co/linoyts/3d_icon_flux_1500#flux-dreambooth-lora---linoyts3d_icon_flux_1500
huggingface.co
linoyts/3d_icon_flux_1500 · Hugging Face
We’re on a journey to advance and democratize artificial intelligence through open source and open science.
Welcome to our interactive lesson on FLUX.1 dev SimpleTuner Test!
I'm excited to help you learn about this topic. Here's a brief overview:
What is FLUX.1 dev SimpleTuner Test?
FLUX.1 dev SimpleTuner Test is a type of artificial intelligence (AI) model that uses a technique called LoRA (Learning to be a Regularizer) to improve the performance of a pre-trained model. In this case, the pre-trained model is called FLUX.1 dev, which is a type of image generation model.
What do you want to learn about FLUX.1 dev SimpleTuner Test?
Here are some options:
What is LoRA and how does it work?
How is FLUX.1 dev SimpleTuner Test different from other image generation models?
What are the key settings and parameters used to train FLUX.1 dev SimpleTuner Test?
How can I use FLUX.1 dev SimpleTuner Test to generate images?
Something else (please specify)
Let me know what you're interested in learning, and I'll do my best to explain it in a way that's easy to understand!
https://huggingface.co/markury/FLUX-dev-LoRA-test
I'm excited to help you learn about this topic. Here's a brief overview:
What is FLUX.1 dev SimpleTuner Test?
FLUX.1 dev SimpleTuner Test is a type of artificial intelligence (AI) model that uses a technique called LoRA (Learning to be a Regularizer) to improve the performance of a pre-trained model. In this case, the pre-trained model is called FLUX.1 dev, which is a type of image generation model.
What do you want to learn about FLUX.1 dev SimpleTuner Test?
Here are some options:
What is LoRA and how does it work?
How is FLUX.1 dev SimpleTuner Test different from other image generation models?
What are the key settings and parameters used to train FLUX.1 dev SimpleTuner Test?
How can I use FLUX.1 dev SimpleTuner Test to generate images?
Something else (please specify)
Let me know what you're interested in learning, and I'll do my best to explain it in a way that's easy to understand!
https://huggingface.co/markury/FLUX-dev-LoRA-test
huggingface.co
markury/FLUX-dev-LoRA-test · Hugging Face
We’re on a journey to advance and democratize artificial intelligence through open source and open science.
Flux DreamBooth LoRA - merve/flux-lego-lora-dreambooth
Model description
These are merve/flux-lego-lora-dreambooth DreamBooth LoRA weights for black-forest-labs/FLUX.1-dev.
The weights were trained using DreamBooth with the Flux diffusers trainer.
Was LoRA for the text encoder enabled? False.
Trigger words
You should use lego set in style of TOK to trigger the image generation.
Download model
Download the *.safetensors LoRA in the Files & versions tab.
Use it with the 🧨 diffusers library https://huggingface.co/merve/flux-lego-lora-dreambooth
Model description
These are merve/flux-lego-lora-dreambooth DreamBooth LoRA weights for black-forest-labs/FLUX.1-dev.
The weights were trained using DreamBooth with the Flux diffusers trainer.
Was LoRA for the text encoder enabled? False.
Trigger words
You should use lego set in style of TOK to trigger the image generation.
Download model
Download the *.safetensors LoRA in the Files & versions tab.
Use it with the 🧨 diffusers library https://huggingface.co/merve/flux-lego-lora-dreambooth
huggingface.co
merve/flux-lego-lora-dreambooth · Hugging Face
We’re on a journey to advance and democratize artificial intelligence through open source and open science.
flux-dreambooth-lora-r16-dev-cfg1
This is a LoRA derived from black-forest-labs/FLUX.1-dev.
The main validation prompt used during training was:
julie, in photograph style
Validation settings
CFG: 3.0
CFG Rescale: 0.0
Steps: 20
Sampler: None
Seed: 420420420
Resolution: 512
Note: The validation settings are not necessarily the same as the training settings.
You can find some example images in the following gallery:https://huggingface.co/ptx0/flux-dreambooth-lora-r16-dev-cfg1
This is a LoRA derived from black-forest-labs/FLUX.1-dev.
The main validation prompt used during training was:
julie, in photograph style
Validation settings
CFG: 3.0
CFG Rescale: 0.0
Steps: 20
Sampler: None
Seed: 420420420
Resolution: 512
Note: The validation settings are not necessarily the same as the training settings.
You can find some example images in the following gallery:https://huggingface.co/ptx0/flux-dreambooth-lora-r16-dev-cfg1
huggingface.co
ptx0/flux-dreambooth-lora-r16-dev-cfg1 · Hugging Face
We’re on a journey to advance and democratize artificial intelligence through open source and open science.
works best when you include 'ps1 game screenshot' in prompt
late 90s/early 2000s ps1/n64 console graphics.
5000 steps
trained on 15 gpt4o captioned and adjusted ps1/n64 game screenshots using https://github.com/ostris/ai-toolkit/tree/main
Trigger words
You should use ps1 to trigger the image generation.
Download model
Weights for this model are available in Safetensors format.
Download them in the Files & versions tab.
Use it with the 🧨 diffusers library
late 90s/early 2000s ps1/n64 console graphics.
5000 steps
trained on 15 gpt4o captioned and adjusted ps1/n64 game screenshots using https://github.com/ostris/ai-toolkit/tree/main
Trigger words
You should use ps1 to trigger the image generation.
Download model
Weights for this model are available in Safetensors format.
Download them in the Files & versions tab.
Use it with the 🧨 diffusers library
flux-dreambooth-lora
This is a LoRA derived from black-forest-labs/FLUX.1-dev.
Two subjects were trained in, a character named Julia (AI) and a real person named River Phoenix.
Empirically, training two subjects in simultaneously kept the model from collapsing, though they don't train evenly - River Phoenix took longer than "Julia", possibly due to the synthetic nature of the data.
The photos of "Julia" came from Flux Pro. River Phoenix images were pulled from Google Image Search, with a focus on high resolution, high quality samples.
No captions were used during training, only instance prompts julia and river phoenix.
The main validation prompt used during training was:
julie, in photograph style
Validation settings
CFG: 3.0
CFG Rescale: 0.0
Steps: 28
Sampler: None
Seed: 420420420
Resolution: 1024x1024
Note: The validation settings are not necessarily the same as the training settings.
You can find some example images in the following gallery:
This is a LoRA derived from black-forest-labs/FLUX.1-dev.
Two subjects were trained in, a character named Julia (AI) and a real person named River Phoenix.
Empirically, training two subjects in simultaneously kept the model from collapsing, though they don't train evenly - River Phoenix took longer than "Julia", possibly due to the synthetic nature of the data.
The photos of "Julia" came from Flux Pro. River Phoenix images were pulled from Google Image Search, with a focus on high resolution, high quality samples.
No captions were used during training, only instance prompts julia and river phoenix.
The main validation prompt used during training was:
julie, in photograph style
Validation settings
CFG: 3.0
CFG Rescale: 0.0
Steps: 28
Sampler: None
Seed: 420420420
Resolution: 1024x1024
Note: The validation settings are not necessarily the same as the training settings.
You can find some example images in the following gallery:
FLUX.1-[dev] Panorama LoRA (v2)
A LoRA model to generate panoramas using Flux dev.
Which image ratio to use?
This model has been trained on images with a 2:1 ratio (2048x1024).
So you might get good results is you use this for width and height.
However panorama viewers are pretty flexible when it comes to resolution, and FLUX.1 seems to generalize well.
For instance the gallery samples have been generated in 1536 × 640 (~21:9), since this is reasonably fast (16 sec on the HF Inference API).
It doesn't work for case X or Y
It usually work fine for "normal" requests, but the model might have trouble creating the spherical distortion if you ask for uncommon content, locations or angles.
If you give me some examples, maybe I can try to find more data to better cover uncommon panoramas.
Non-commercial use
As the base model is FLUX.1-[dev] and since the data comes from Google Street View, it should be used for non-commercial, personal or demonstration purposes only.
Please use it responsibly, thank you!
A LoRA model to generate panoramas using Flux dev.
Which image ratio to use?
This model has been trained on images with a 2:1 ratio (2048x1024).
So you might get good results is you use this for width and height.
However panorama viewers are pretty flexible when it comes to resolution, and FLUX.1 seems to generalize well.
For instance the gallery samples have been generated in 1536 × 640 (~21:9), since this is reasonably fast (16 sec on the HF Inference API).
It doesn't work for case X or Y
It usually work fine for "normal" requests, but the model might have trouble creating the spherical distortion if you ask for uncommon content, locations or angles.
If you give me some examples, maybe I can try to find more data to better cover uncommon panoramas.
Non-commercial use
As the base model is FLUX.1-[dev] and since the data comes from Google Street View, it should be used for non-commercial, personal or demonstration purposes only.
Please use it responsibly, thank you!
Here is my first crack at a realism Flux model, and I guess the first realism model I've ever shared publicly.
All the training data was opensource/open license photographs. I do intend to revisit it and expand and improve on it. It can benefit from emphasizing the style by adding "vintage" or "faded film" to the prompt.
Big appreciation to Glif for sponsoring the model!
All the training data was opensource/open license photographs. I do intend to revisit it and expand and improve on it. It can benefit from emphasizing the style by adding "vintage" or "faded film" to the prompt.
Big appreciation to Glif for sponsoring the model!
Flux Dev 1 model that creates images with both photographic and illustrative elements. Pretty cool right? I've trained this model on a curated collection of images gathered from Pinterest. I am not the creator of the original art used for the fine-tuned training, but want to provide access for creative exploration. You can run the model for yourself or train your own on Replicate. Let me know your thoughts and have a wonderful day!
https://replicate.com/lucataco/ai-toolkit/readme
https://replicate.com/lucataco/flux-dev-lora
Trigger words
It would help if you used in the style of TOK for better style preservation. It is best to place the trigger words first and describe illustrative elements in your scene like clothing or expressive elements.
Download model
Download the *.safetensors LoRA in the Files & versions tab.
Use it with the 🧨 diffusers library
https://replicate.com/lucataco/ai-toolkit/readme
https://replicate.com/lucataco/flux-dev-lora
Trigger words
It would help if you used in the style of TOK for better style preservation. It is best to place the trigger words first and describe illustrative elements in your scene like clothing or expressive elements.
Download model
Download the *.safetensors LoRA in the Files & versions tab.
Use it with the 🧨 diffusers library
Replicate
lucataco/ai-toolkit – Replicate
Ostris AI-Toolkit for Flux LoRA Training
Stable-Diffusion: Optimized for Mobile Deployment
State-of-the-art generative AI model used to generate detailed images conditioned on text descriptions
Generates high resolution images from text prompts using a latent diffusion model. This model uses CLIP ViT-L/14 as text encoder, U-Net based latent denoising, and VAE based decoder to generate the final image.
This model is an implementation of Stable-Diffusion found here. This repository provides scripts to run Stable-Diffusion on Qualcomm® devices. More details on model performance across various devices, can be found here.https://huggingface.co/qualcomm/Stable-Diffusion
State-of-the-art generative AI model used to generate detailed images conditioned on text descriptions
Generates high resolution images from text prompts using a latent diffusion model. This model uses CLIP ViT-L/14 as text encoder, U-Net based latent denoising, and VAE based decoder to generate the final image.
This model is an implementation of Stable-Diffusion found here. This repository provides scripts to run Stable-Diffusion on Qualcomm® devices. More details on model performance across various devices, can be found here.https://huggingface.co/qualcomm/Stable-Diffusion
huggingface.co
qualcomm/Stable-Diffusion · Hugging Face
We’re on a journey to advance and democratize artificial intelligence through open source and open science.
Score-Based Generative Modeling through Stochastic Differential Equations (SDE)
Paper: Score-Based Generative Modeling through Stochastic Differential Equations
Authors: Yang Song, Jascha Sohl-Dickstein, Diederik P. Kingma, Abhishek Kumar, Stefano Ermon, Ben Poole
Abstract:
Creating noise from data is easy; creating data from noise is generative modeling. We present a stochastic differential equation (SDE) that smoothly transforms a complex data distribution to a known prior distribution by slowly injecting noise, and a corresponding reverse-time SDE that transforms the prior distribution back into the data distribution by slowly removing the noise. Crucially, the reverse-time SDE depends only on the time-dependent gradient field (\aka, score) of the perturbed data distribution. By leveraging advances in score-based generative modeling, we can accurately estimate these scores with neural networks, and use numerical SDE solvers to generate samples. We show that this framework encapsulates previous approaches in score-based generative modeling and diffusion probabilistic modeling, allowing for new sampling procedures and new modeling capabilities. In particular, we introduce a predictor-corrector framework to correct errors in the evolution of the discretized reverse-time SDE. We also derive an equivalent neural ODE that samples from the same distribution as the SDE, but additionally enables exact likelihood computation, and improved sampling efficiency. In addition, we provide a new way to solve inverse problems with score-based models, as demonstrated with experiments on class-conditional generation, image inpainting, and colorization. Combined with multiple architectural improvements, we achieve record-breaking performance for unconditional image generation on CIFAR-10 with an Inception score of 9.89 and FID of 2.20, a competitive likelihood of 2.99 bits/dim, and demonstrate high fidelity generation of 1024 x 1024 images for the first time from a score-based generative model.
Paper: Score-Based Generative Modeling through Stochastic Differential Equations
Authors: Yang Song, Jascha Sohl-Dickstein, Diederik P. Kingma, Abhishek Kumar, Stefano Ermon, Ben Poole
Abstract:
Creating noise from data is easy; creating data from noise is generative modeling. We present a stochastic differential equation (SDE) that smoothly transforms a complex data distribution to a known prior distribution by slowly injecting noise, and a corresponding reverse-time SDE that transforms the prior distribution back into the data distribution by slowly removing the noise. Crucially, the reverse-time SDE depends only on the time-dependent gradient field (\aka, score) of the perturbed data distribution. By leveraging advances in score-based generative modeling, we can accurately estimate these scores with neural networks, and use numerical SDE solvers to generate samples. We show that this framework encapsulates previous approaches in score-based generative modeling and diffusion probabilistic modeling, allowing for new sampling procedures and new modeling capabilities. In particular, we introduce a predictor-corrector framework to correct errors in the evolution of the discretized reverse-time SDE. We also derive an equivalent neural ODE that samples from the same distribution as the SDE, but additionally enables exact likelihood computation, and improved sampling efficiency. In addition, we provide a new way to solve inverse problems with score-based models, as demonstrated with experiments on class-conditional generation, image inpainting, and colorization. Combined with multiple architectural improvements, we achieve record-breaking performance for unconditional image generation on CIFAR-10 with an Inception score of 9.89 and FID of 2.20, a competitive likelihood of 2.99 bits/dim, and demonstrate high fidelity generation of 1024 x 1024 images for the first time from a score-based generative model.